5 released by RunwayML. Take the “Behind the scenes of the moon landing” image. Check it out: Stable Diffusion Photoshop Plugin (0. Download any of the VAEs listed above and place them in the folder stable-diffusion-webuimodelsVAE. #. You'll see this on the txt2img tab:You can make NSFW images In Stable Diffusion using Google Colab Pro or Plus. Syntax: cv2. Running App Files Files Community 37. マイクロソフトは DirectML を最適化し、Stable Diffusion で使用されているトランスフォーマーと拡散モデルを高速化することで、Windows ハードウェア・エコシステム全体でより優れた動作を実現しました。 AMD は、Olive のプレリリースに見られるように. Steps. Set sampling steps to 20 and sampling method to DPM++ 2M Karras. ckpt (5. Goodbye Babel, generated by Andrew Zhu using Diffusers in pure Python. ckpt) Place the model file inside the modelsstable-diffusion directory of your installation directory (e. txt2img2img is an. Also there is post tagged here where all the links to all resources are. Pak jsem si řekl, že zkusím img2txt a ten vytvořil. Want to see examples of what you can build with Replicate? Check out our showcase. These are our findings: Many consumer grade GPUs can do a fine job, since stable diffusion only needs about 5 seconds and 5 GB of VRAM to run. 24, so if you have that or a newer version, you don't need the workaround anymore. To differentiate what task you want to use the checkpoint for, you have to load it directly with its corresponding task-specific pipeline class:La manera más sencilla de utilizar Stable Diffusion es registrarte en un editor de imágenes por IA llamado Dream Studio. stablediffusiononw. information gathering ; txt2img ; img2txt ; stable diffusion ; Stable Diffusion is a tool to create pictures with keywords. By my understanding, a lower value will be more "creative" whereas a higher value will adhere more to the prompt. ” img2img ” diffusion) can be a powerful technique for creating AI art. Features. openai. DreamBooth. All you need to do is to use img2img method, supply a prompt, dial up the CFG scale, and tweak the denoising strength. Settings for all eight stayed the same: Steps: 20, Sampler: Euler a, CFG scale: 7, Face restoration: CodeFormer, Size: 512x768, Model hash: 7460a6fa. stable-diffusion-img2img. Notice there are cases where the output is barely recognizable as a rabbit. 5] Since, I am using 20 sampling steps, what this means is using the as the negative prompt in steps 1 – 10, and (ear:1. /webui. Number of denoising steps. File "C:UsersGros2stable-diffusion-webuildmmodelslip. Useful resource. In this quick episode we do a simple workflow where we upload an image into our SDXL graph inside of ComfyUI and add additional noise to produce an altered i. ControlNet is a brand new neural network structure that allows, via the use of different special models, to create image maps from any images and using these. Generate the image. 5. Stable Diffusion is a concealed text-to-image diffusion model, capable of generating photorealistic images from any textual input, fosters independent flexibility in producing remarkable visuals. C:stable-diffusion-uimodelsstable-diffusion)Option 1: Every time you generate an image, this text block is generated below your image. That’s the basic. g. 5를 그대로 사용하며, img2txt. There is no rule here - the more area of the original image is covered, the better match. txt2img, img2img, depth2img, pix2pix, inpaint and interrogation (img2txt). Don't use other versions unless you are looking for trouble. Intro to ComfyUI. Predictions typically complete within 14 seconds. A checker for NSFW images. Resize and fill: This will add in new noise to pad your image to 512x512, then scale to 1024x1024, with the expectation that img2img will. Download and install the latest Git here. 📚 RESOURCES- Stable Diffusion web de. Render: the act of transforming an abstract representation of an image into a final image. 因為是透過 Stable Diffusion Model 算圖，除了放大解析度外，還能增加細部細節！. Please reopen this issue! Deleting config. stable-diffusion txt2img参数整理 Sampling steps ：采样步骤”：“迭代改进生成图像的次数；较高的值需要更长的时间；非常低的值可能会产生糟糕的结果”， 指的是Stable Diffusion生成图像所需的迭代步数。Stable Diffusion is a cutting-edge text-to-image diffusion model that can generate photo-realistic images based on any given text input. It generates accurate, diverse and creative captions for images. I am still new to Stable Diffusion, but I still managed to get an art piece with text, nonetheless. You can also upload and replicate non-AI generated images. Transform your doodles into real images in seconds. It is an effective and efficient approach that can be applied to image understanding in numerous scenarios, especially when examples are scarce. 160 upvotes · 39 comments. The script outputs an image file based on the model's interpretation of the prompt. If you don't like the results, you can generate new designs an infinite number of times until you find a logo you absolutely love! Watch It In Action. . I have been using Stable Diffusion for about 2 weeks now. For the rest of this guide, we'll either use the generic Stable Diffusion v1. The same issue occurs if an image with a variation seed is created on the txt2img tab and the "Send to img2txt" option is used. 使用管理员权限打开下图应用程序. Pipeline for text-to-image generation using Stable Diffusion. But the […]Stable Diffusion是2022年發布的深度學習 文本到图像生成模型。 它主要用於根據文本的描述產生詳細圖像，儘管它也可以應用於其他任務，如內補繪製、外補繪製，以及在提示詞指導下產生圖生圖的转变。. Hieronymus Bosch. A text-to-image generative AI model that creates beautiful images. 가장먼저 파이썬이라는 프로그램이 돌아갈 수 있도록 Python을 설치합니다. /. Python. g. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. ckpt files) must be separately downloaded and are required to run Stable Diffusion. 4 ・diffusers 0. Put the Lora of the first epoch in your prompt (like "<lora:projectname-01:0. SFW and NSFW generations. 1. Sort of new here. Stable diffusion image-to-text (SDIT) is an advanced image captioning model based on the GPT architecture and uses a diffusion-based training algorithm to improve stability and consistency during training. Mage Space and Yodayo are my recommendations if you want apps with more social features. 9% — contains NSFW material, giving the model little to go on when it comes to explicit content. Files to download:👉Python: dont have the stable-diffusion-v1 folder, i have a bunch of others tho. 手順3：学習を行う. It serves as a quick reference as to what the artist's style yields. To shed light on these questions, we present an inference benchmark of Stable Diffusion on different GPUs and CPUs. At least that is what he says. Appendix A: Stable Diffusion Prompt Guide. Text to image generation. TurbTastic •. When it comes to speed to output a single image, the most powerful. Discover amazing ML apps made by the communityPosition the 'Generation Frame' in the right place. com) r/StableDiffusion. py", line 144, in interrogate load_blip_model(). jpeg by default on the root of the repo. zip. How to use ChatGPT. Press the Window key (It should be on the left of the space bar on your keyboard), and a search window should appear. 部署 Stable Diffusion WebUI . Textual inversion is NOT img2txt! Let's make sure people don't start calling img2txt textual inversion, because these things are two completely different applications. A buddy of mine told me about it being able to be locally installed on a machine. Aug 26, 2022. You can receive up to four options per prompt. Press Send to img2img to send this image and parameters for outpainting. Unprompted is a highly modular extension for AUTOMATIC1111's Stable Diffusion Web UI that allows you to include various shortcodes in your prompts. Additionally, their formulation allows to apply them to image modification tasks such as inpainting directly without retraining. At the time of release (October 2022), it was a massive improvement over other anime models. For certain inputs, simply running the model in a convolutional fashion on larger features than it was trained on can sometimes result in interesting results. It is common to use negative embeddings for anime. 7>"), and on the script's X value write something like "-01, -02, -03", etc. card. creates original designs within seconds. It allows the model to generate contextualized images of the subject in different scenes, poses, and views. Another experimental VAE made using the Blessed script. London- and California-based startup Stability AI has released Stable Diffusion, an image-generating AI that can produce high-quality images that look as if they were. The Payload config is central to everything that Payload does. It came out gibberish though. LoRAモデルを使って画像を生成する方法（Stable Diffusion web UIが必要）. Troubleshooting. Compress the prompt and fixes. All the training scripts for text-to-image finetuning used in this guide can be found in this repository if you’re interested in taking a closer look. The tool then processes the image using its stable diffusion algorithm and generates the corresponding text output. If you’ve saved new models in there while A1111 is running you can hit the blue refresh button to the right of the drop. Rising. Negative embeddings bad artist and bad prompt. and i'll got a same problem again and again Stable diffusion model failed to load, exiting. Dear friends, come and join me on an incredible journey through Stable Diffusion. The extensive list of features it offers can be intimidating. 打开stable-diffusion-webuimodelsstable-diffusion目录，此处为各种模型的存放处。 需要预先存放一个模型才能正常使用。 3. 5. The last model containing NSFW concepts was 1. run. Payload is a config-based, code-first CMS and application framework. I. Stable Diffusion v1. In this tutorial I’ll cover: A few ways this technique can be useful in practice. Run time and cost. 指定した画像に近づくように画像生成する機能です。通常のプロンプトによる生成指定に加えて、追加でVGG16の特徴量を取得し、生成中の画像が指定したガイド画像に近づくよう、生成される画像をコントロールします。 2. Navigate to txt2img tab, find Amazon SageMaker Inference panel. Wait a few moments, and you'll have four AI-generated options to choose from. The maximum value is 4. ，AI绘画stable diffusion，AI辅助室内设计controlnet-语义分割控制测试-3. Img2Prompt. See the SDXL guide for an alternative setup with SD. The model bridges the gap between vision and natural. To put another way, quoting your source at gigazine, "the larger the CFG scale, the more likely it is that a new image can be generated according to the image input by the prompt. Introducing Stable Fast: An ultra lightweight inference optimization library for HuggingFace Diffusers on NVIDIA GPUs r/linuxquestions • How to install gcc-arm-linux-gnueabihf 4. It’s trained on 512x512 images from a subset of the LAION-5B dataset. If the image with the text was clear enough, you will receive recognized and readable text. This version is optimized for 8gb of VRAM. 20. This process is called "reverse diffusion," based on math inspired. with current technology would it be possible to ask the AI to generate a text from an image? in order to know what technology could describe the image, a tool for AI to describe the image for us. The backbone. Step 3: Clone web-ui. Roboti na kole. This guide will show you how to finetune the CompVis/stable-diffusion-v1-4 model on your own dataset with PyTorch and Flax. Improving image generation at different aspect ratios using conditional masking during training. The original Stable Diffusion model was created in a collaboration with CompVis and RunwayML and builds upon the work: High-Resolution Image Synthesis with Latent Diffusion Models. Just two. To use a VAE in AUTOMATIC1111 GUI, go to the Settings tab and click the Stabe Diffusion section on the left. . (Optimized for stable-diffusion (clip ViT-L/14)) Public; 2. OCR or Optical Character Recognition has never been so easy. Cung cấp bộ công cụ và hướng dẫn hoàn toàn miễn phí, giúp bất kỳ cá nhân nào cũng có thể tiếp cận được công cụ vẽ tranh AI Stable DiffusionFree Stable Diffusion webui - txt2img img2img. Sort of new here. Create multiple variants of an image with Stable Diffusion. By default, 🤗 Diffusers automatically loads these . Negative prompting influences the generation process by acting as a high-dimension anchor,. The Stable Diffusion model was created by researchers and engineers from CompVis, Stability AI, Runway, and LAION. 5. To use this pipeline for image-to-image, you’ll need to prepare an initial image to pass to the pipeline. Share Tweak it. . A decoder, which turns the final 64x64 latent patch into a higher-resolution 512x512 image. Stable Diffusion 2. ; Download the optimized Stable Diffusion project here. ; Mind you, the file is over 8GB so while you wait for the download. Stable Diffusionのプロンプトは英文に近いものですので、作成をChatGPTに任せることは難しくないはずです。. be 131 upvotes · 15 comments StableDiffusion. However, there’s a twist. Steps. k. By simply replacing all instances linking to the original script with the script that has no safety filters, you can easily achieve generate NSFW images. Explore and run machine. Items you don't want in the image. 9) in steps 11-20. You'll have a much easier time if you generate the base image in SD, add in text with a conventional image editing program. The default value is set to 2. We would like to show you a description here but the site won’t allow us. Hey there! I’ve been doing some extensive tests between diffuser’s stable diffusion and AUTOMATIC1111’s and NMKD-SD-GUI implementations (which both wrap the CompVis/stable-diffusion repo). An image generated at resolution 512x512 then upscaled to 1024x1024 with Waifu Diffusion 1. Spaces. 08:08. Predictions typically complete within 2 seconds. Number of images to be returned in response. The VD-basic is an image variation model with a single-flow. The layout of Stable Diffusion in DreamStudio is more cluttered than DALL-E 2 and Midjourney, but it's still easy to use. I. This model uses a frozen CLIP ViT-L/14 text. ·. 购买云端服务器-> 内网穿透 -> api形式运行sd -> 手机发送api请求，即可实现. Start the WebUI. The average face of a teacher generated by Stable Diffusion and DALL-E 2. I've been using it to add pictures to any of the recipes that are added to my wiki site without a picture. Functioning as image viewers for the terminal, chafa and catimg have only been an integral part of a stable release of the Linux distribution since Debian GNU/Linux 10. Find your API token in your account settings. 本视频基于AI绘图软件Stable Diffusion。. Inpainting appears in the img2img tab as a seperate sub-tab. Txt2Img：文生图 Img2Txt：图生文 Img2Img：图生图 功能点 部署 Stable Diffusion WebUI 更新 python 版本 切换国内 Linux 安装镜像 安装 Nvidia 驱动 安装stable-diffusion-webui 并启动服务 部署飞书机器人 操作方式 操作命令 设置关键词: 探索企联AI Hypernetworks. 81 seconds. Para ello vam. Creating applications on Stable Diffusion’s open-source platform has proved wildly successful. Commit where the problem happens. Contents. Using stable diffusion and these prompts hand-in-hand, you can easily create stunning and high-quality logos in seconds without needing any design experience. 9 fine, but when I try to add in the stable-diffusion. CLIP Interrogator extension for Stable Diffusion WebUI. 0 (SDXL 1. Mage Space has very limited free features, so it may as well be a paid app. I found a genius who uses ControlNet and OpenPose to change the poses of pixel art character! self. 5 it/s. stable diffusion webui 脚本使用方法（下），人脸编辑还不错. Check it out: Stable Diffusion Photoshop Plugin (0. 1M runs. Go to Settings tab. card classic compact. Img2txt. Go to img2txt tab. Those are the absolute minimum system requirements for Stable Diffusion. 이제 부터 Stable Diffusion은 줄여서 SD로 표기하겠습니다. The StableDiffusionPipeline is capable of generating photorealistic images given any text input. But the width, height and other defaults need changing. A Keras / Tensorflow implementation of Stable Diffusion. MORPH_CLOSE, kernel) -> image: Input Image array. exe, follow instructions. This is a repo providing same stable diffusion experiments, regarding textual inversion task and captioning task pytorch clip captioning-images img2txt caption-generation caption-generator huggingface latent-diffusion stable-diffusion huggingface-diffusers latent-diffusion-models textual-inversion VGG16 Guided Stable Diffusion. img2img settings. 4 min read. On SD 2. ago. Preview. r/StableDiffusion •. Put this in the prompt text box. env. SDXL is a larger and more powerful version of Stable Diffusion v1. Next, VD-DC is a two-flow model that supports both text-to-image synthesis and image-variation. This parameter controls the number of these denoising steps. Set the batch size to 4 so that you can. Check out the img2img. You can use them to remove specific elements, styles, or. Hosted on Banana 🍌. Live Demo at Available on Hugging Facesuccinctly/text2image-prompt-generatorlike229. jkcarney commented Jun 30, 2023. BLIP: image used in this demo is from Stephen Young: #3: Using Stable Diffusion’s PNG Info. 5. To use this, first make sure you are on latest commit with git pull, then use the following command line argument: In the img2img tab, a new button will be available saying "Interrogate DeepBooru", drop an image in and click the button. Copy the prompt, paste it to the Stable Diffusion and press Generate to see generated images. Search Results related to img2txt. ChatGPT is aware of the history of your current conversation. It is a parameter that tells the Stable Diffusion model what not to include in the generated image. 丨Stable Diffusion终极教程【第5期】，Stable Diffusion提示词起手式TAG（中文界面），DragGAN真有那么神？在线运行 + 开箱评测。，Stable Diffusion教程之animatediff生成丝滑动画（一），【简易化】finetune定制大模型， Dreambooth webui画风训练保姆教程，当ai水说话开始喘气. 6. 除了告訴 Stable Diffusion 有哪些物品，亦可多加該物的形容詞，如人的穿著、動作、年齡等等描述; 地：物體所在地，亦可想像成畫面的背景，讓 Stable Diffusion 知道背景要畫什麼（不然他會自由發揮） 風格：告訴 Stable Diffusion 要以什麼風格呈現圖片，某個畫家？ Stable Diffusion WebUI (AUTOMATIC1111 or A1111 for short) is the de facto GUI for advanced users. I used two different yet similar prompts and did 4 A/B studies with each prompt. It’s a simple and straightforward process that doesn’t require any technical expertise. Linux: run the command webui-user. . (Open in Colab) Build your own Stable Diffusion UNet model from scratch in a notebook. Dreamshaper. 04 and probably any later versions with ImageMagick 6, here's how you fix the issue by removing that workaround:. Hot New Top Rising. Take careful note of the syntax of the example that’s already there. In this step-by-step tutorial, learn how to download and run Stable Diffusion to generate images from text descriptions. a. Note: This repo aims to provide a Ready-to-Go setup with TensorFlow environment for Image Captioning Inference using pre-trained model. for examples:"Logo of a pirate","logo of a sunglass with girl" or something complex like "logo of a ice-cream with snake" etc. Let's dive in deep and learn how to generate beautiful AI Art based on prom. Stable Diffusion Prompts Generator helps you. 16:17. The latest stability ai release is 2. In closing operation, the basic premise is that the closing is opening performed in reverse. com uses a Commercial suffix and it's server(s) are located in N/A with the IP number 104. AUTOMATIC1111のモデルデータは「"stable-diffusion-webuimodelsStable-diffusion"」の中にあります。 正則化画像の用意. Note: Earlier guides will say your VAE filename has to have the same as your model filename. Updated 1 day, 17 hours ago 53 runs fofr / sdxl-pixar-cars SDXL fine-tuned on Pixar Cars. 使用anaconda进行webui的创建. Stable Diffusion img2img support comes to Photoshop. ckpt Global Step: 140000 Traceback (most recent call last): File "D:AIArtstable-diffusion-webuivenvlibsite. true. ai says it can double the resolution of a typical 512×512 pixel image in half a second. It’s trained on 512x512 images from a subset of the LAION-5B dataset. 4M runs. ControlNet is a neural network structure to control diffusion models by adding extra conditions. . 🙏 Thanks JeLuF for providing these directions. Prompt: the description of the image the AI is going to generate. Save a named theme "Chris's 768". ckpt for using v1. As of June 2023, Midjourney also gained inpainting and outpainting via the Zoom Out button. Subsequently, to relaunch the script, first activate the Anaconda command window (step 3), enter the stable-diffusion directory (step 5, "cd path ostable-diffusion"), run "conda activate ldm" (step 6b), and then launch the dream script (step 9). About that huge long negative prompt list. 5 anime-like image generations. Stability AI는 방글라데시계 영국인. DiffusionBee is one of the easiest ways to run Stable Diffusion on Mac. The CLIP interrogator has two parts: one is the BLIP model, which takes on the function of decoding and reasoning about the text description. If there is a text-to-image model that can come very close to Midjourney, then it’s Stable Diffusion. You can also upload and replicate non-AI generated images. Additional training is achieved by training a base model with an additional dataset you are. 5 model. 比如我的路径是D:dataicodinggit_hubdhumanstable-diffusion-webuimodelsStable-diffusion 在项目目录内安装虚拟环境 python -m venv venv_port 执行webui-user. BLIP-2 is a zero-shot visual-language model that can be used for multiple image-to-text tasks with image and image and text prompts. 26. The easiest way to try it out is to use one of the Colab notebooks: ; GPU Colab ; GPU Colab Img2Img ; GPU Colab Inpainting ; GPU Colab - Tile / Texture generation ; GPU Colab - Loading. txt2img OR "imaging" is mathematically divergent operation, from less bits to more bits, even ARM or RISC-V can do that. Run time and cost. More awesome work from Christian Cantrell in his free plugin. AI画像生成士. Example outputs . 0 was released in November 2022 and has been entirely funded and developed by Stability AI. Get an approximate text prompt, with style, matching an image. Doing this on a loop takes advantage of the imprecision in using CLIP latent space walk - fixed seed but two different prompts. ago. You can use the. . The company claims this is the fastest-ever local deployment of the tool on a smartphone. Enjoy . 1:7860" or "localhost:7860" into the address bar, and hit Enter. If you are using any of the popular WebUI stable diffusions (like Automatic1111) you can use Inpainting. Installing. Using a model is an easy way to achieve a certain style. Stable Diffusion - Image to Prompts Run 934. 10. StabilityAI’s Stable Video Diffusion (SVD), image to video Updated 4 hours ago 173 runs sdxl A text-to-image generative AI model that creates beautiful images Updated 2 weeks, 2 days ago 20. (com a tecnologia atual seria possivel solicitar a IA gerar um texto a partir de uma imagem ? com a finalidade de saber o que a tecnologia poderia. Stable Diffusion without UI or tricks (only take off filter xD). Stable Diffusion web UIをインストールして使えるようにしておく。 Stable Diffusion web UI用のControlNet拡張機能もインストールしておく。 この2つについては下記の記事でやり方等を丁寧にご説明していますので、まだ準備ができていないよという方はそちらも併せて. ckpt file was a choice. What is Img2Img in Stable Diffusion Setting up The Software for Stable Diffusion Img2img How to Use img2img in Stable Diffusion Step 1: Set the background Step 2: Draw the Image Step 3: Apply Img2Img The End! For those who haven’t been blessed with innate artistic abilities, fear not! Img2Img and Stable Diffusion can. On the first run, the WebUI will download and install some additional modules. ネットにあるあの画像、私も作りたいな〜. Hosted on Banana 🍌. A method to fine tune weights for CLIP and Unet, the language model and the actual image de-noiser used by Stable Diffusion, generously donated to the world by our friends at Novel AI in autumn 2022. The image and prompt should appear in the img2img sub-tab of the img2img tab. 2. py file for more options, including the number of steps. Images generated by Stable Diffusion based on the prompt we’ve. 002. 1 I use this = oversaturated, ugly, 3d, render, cartoon, grain, low-res, kitsch, black and white. 手順2：「gui. The inspiration was simply the lack of any Emiru model of any sort here. Start with installation & basics, then explore advanced techniques to become an expert. More awesome work from Christian Cantrell in his free plugin. (You can also experiment with other models. You can create your own model with a unique style if you want. Also, because the Payload source code is fully written in. Then you can either mask the face and choose inpaint unmasked, or select only the parts you want changed and inpaint masked. Learn the importance, workings, and benefits of using Kiwi Prompt's chat GPT & Google Bard prompts to enhance your stable diffusion writing. With those sorts of specs, you. Textual Inversion is a technique for capturing novel concepts from a small number of example images. Prompt by Rachey13x 17 days ago (8k, RAW photo, highest quality), hyperrealistic, Photo of a gang member from Peaky Blinders on a hazy and smokey dark alley, highly detailed, cinematic, film. - use img2txt to generate the prompt and img2img to provide the starting point. Hiresは「High Resolution」の略称で高解像度という意味を持ち、fixは「修正・変更」を指します。. I have showed you how easy it is to use Stable Diffusion to stylize images. Select. Download: Installation: Extract anywhere (not a protected folder - NOT Program Files - preferrably a short custom path like D:/Apps/AI/), run StableDiffusionGui. Stable Difussion Web UIのHires. stable diffusion webui 脚本使用方法（上）. AUTOMATIC1111 Web-UI is a free and popular Stable Diffusion software. 4 (v1. Discover amazing ML apps made by the communitystability-ai / stable-diffusion. 0-base. From left to right, top to bottom: Lady Gaga, Boris Johnson, Vladimir Putin, Angela Merkel, Donald Trump, Plato. I was using one but it does not work anymore since yesterday. Flirty_Dane • 7 mo. Having the Stable Diffusion model and even Automatic’s Web UI available as open-source is an important step to democratising access to state-of-the-art AI tools. Stability. Easy Prompt SelectorのYAMLファイルは「stable-diffusion-webuiextensionssdweb-easy-prompt-selector ags」の中にあります。 「. Documentation is lacking. Search. 0. . Uncrop. Is there an alternative. img2txt arch. ai and more. For training from scratch or funetuning, please refer to Tensorflow Model Repo. ps1」を実行して設定を行う. PromptMateIO • 7 mo. Stable Diffusion consists of three parts: A text encoder, which turns your prompt into a latent vector. If you put your picture in, would Stable Diffusion start roasting you with tags?.